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ComfyUI Text To Image Workflow

This document provides a comprehensive guide on using ComfyUI for text-to-image AI art generation, outlining the workflow, necessary components, and principles of diffusion models. It details the steps to set up the workflow, load models, and generate images, along with prompting techniques for optimal results. Additionally, it introduces the SD1.5 model, highlighting its advantages and limitations in the context of AI art generation.

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0% found this document useful (0 votes)
173 views15 pages

ComfyUI Text To Image Workflow

This document provides a comprehensive guide on using ComfyUI for text-to-image AI art generation, outlining the workflow, necessary components, and principles of diffusion models. It details the steps to set up the workflow, load models, and generate images, along with prompting techniques for optimal results. Additionally, it introduces the SD1.5 model, highlighting its advantages and limitations in the context of AI art generation.

Uploaded by

sculp4all
Copyright
© © All Rights Reserved
We take content rights seriously. If you suspect this is your content, claim it here.
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Download as PDF, TXT or read online on Scribd
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ComfyUI Text to Image Workflow

This guide will help you understand the concept of text-to-image in AI art generation and
complete a text-to-image workflow in ComfyUI

This guide aims to introduce you to ComfyUI's text-to-image workflow and help you understand
the functionality and usage of various ComfyUI nodes.

In this document, we will:

Complete a text-to-image workflow


Gain a basic understanding of diffusion model principles
Learn about the functions and roles of workflow nodes
Get an initial understanding of the SD1.5 model

We'll start by running a text-to-image workflow, followed by explanations of related concepts.


Please choose the relevant sections based on your needs.

About Text to Image


Text to Image is a fundamental process in AI art generation that creates images from text
descriptions, with diffusion models at its core.

The text-to-image process requires the following elements:

Artist: The image generation model


Canvas: The latent space
Image Requirements (Prompts): Including positive prompts (elements you want in the
image) and negative prompts (elements you don't want)

This text-to-image generation process can be simply understood as telling your requirements
(positive and negative prompts) to an artist (the image model), who then creates what you
want based on these requirements.

ComfyUI Text to Image Workflow Example Guide


1. Preparation
Ensure you have at least one SD1.5 model file in your ComfyUI/models/checkpoints folder,
such as v1-5-pruned-emaonly-fp16.safetensors
If you haven't installed it yet, please refer to the model installation section in Getting Started
with ComfyUI AI Art Generation.

2. Loading the Text to Image Workflow


Download the image below and drag it into ComfyUI to load the workflow:

3. Loading the Model and Generating Your First Image


After installing the image model, follow the steps in the image below to load the model and
generate your first image
Follow these steps according to the image numbers:

1. In the Load Checkpoint node, use the arrows or click the text area to ensure v1-5-pruned-
emaonly-fp16.safetensors is selected, and the left/right arrows don't show null text
2. Click the Queue button or use the shortcut Ctrl + Enter to execute image generation

After the process completes, you should see the resulting image in the Save Image node
interface, which you can right-click to save locally
If you're not satisfied with the result, try running the generation multiple times. Each time you
run it, KSampler will use a different random seed based on the seed parameter, so each
generation will produce different results

4. Start Experimenting
Try modifying the text in the CLIP Text Encoder
The Positive connection to the KSampler node represents positive prompts, while the
Negative connection represents negative prompts

Here are some basic prompting principles for the SD1.5 model:

Use English whenever possible


Separate prompts with English commas ,
Use phrases rather than long sentences
Use specific descriptions
Use expressions like (golden hour:1.2) to increase the weight of specific keywords,
making them more likely to appear in the image. 1.2 is the weight, golden hour is the
keyword
Use keywords like masterpiece, best quality, 4k to improve generation quality

Here are several prompt examples you can try, or use your own prompts for generation:

1. Anime Style

Positive prompts:
anime style, 1girl with long pink hair, cherry blossom background, studio
ghibli aesthetic, soft lighting, intricate details

masterpiece, best quality, 4k

Negative prompts:

low quality, blurry, deformed hands, extra fingers

2. Realistic Style

Positive prompts:

(ultra realistic portrait:1.3), (elegant woman in crimson silk dress:1.2),


full body, soft cinematic lighting, (golden hour:1.2),
(fujifilm XT4:1.1), shallow depth of field,
(skin texture details:1.3), (film grain:1.1),
gentle wind flow, warm color grading, (perfect facial symmetry:1.3)

Negative prompts:

(deformed, cartoon, anime, doll, plastic skin, overexposed, blurry, extra


fingers)

3. Specific Artist Style

Positive prompts:

fantasy elf, detailed character, glowing magic, vibrant colors, long flowing
hair, elegant armor, ethereal beauty, mystical forest, magical aura, high
detail, soft lighting, fantasy portrait, Artgerm style

Negative prompts:

blurry, low detail, cartoonish, unrealistic anatomy, out of focus, cluttered,


flat lighting

Text to Image Working Principles


The entire text-to-image process can be understood as a reverse diffusion process. The v1-
5-pruned-emaonly-fp16.safetensors we downloaded is a pre-trained model that can generate
target images from pure Gaussian noise. We only need to input our prompts, and it can
generate target images through denoising random noise.

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We need to understand two concepts:

1. Latent Space: Latent Space is an abstract data representation method in diffusion models.
Converting images from pixel space to latent space reduces storage space and makes it
easier to train diffusion models and reduce denoising complexity. It's like architects using
blueprints (latent space) for design rather than designing directly on the building (pixel
space), maintaining structural features while significantly reducing modification costs
2. Pixel Space: Pixel Space is the storage space for images, which is the final image we see,
used to store pixel values.

If you want to learn more about diffusion models, you can read these papers:

Denoising Diffusion Probabilistic Models (DDPM)


Denoising Diffusion Implicit Models (DDIM)
High-Resolution Image Synthesis with Latent Diffusion Models

ComfyUI Text to Image Workflow Node Explanation


A. Load Checkpoint Node

This node is typically used to load the image generation model. A checkpoint usually contains
three components: MODEL (UNet) , CLIP , and VAE

MODEL (UNet) : The UNet model responsible for noise prediction and image generation
during the diffusion process
CLIP : The text encoder that converts our text prompts into vectors that the model can
understand, as the model cannot directly understand text prompts
VAE : The Variational AutoEncoder that converts images between pixel space and latent
space, as diffusion models work in latent space while our images are in pixel space

B. Empty Latent Image Node

Defines a latent space that outputs to the KSampler node. The Empty Latent Image node
constructs a pure noise latent space

You can think of its function as defining the canvas size, which determines the dimensions of
our final generated image

C. CLIP Text Encoder Node


Used to encode prompts, which are your requirements for the image

The Positive condition input connected to the KSampler node represents positive
prompts (elements you want in the image)
The Negative condition input connected to the KSampler node represents negative
prompts (elements you don't want in the image)

The prompts are encoded into semantic vectors by the CLIP component from the Load
Checkpoint node and output as conditions to the KSampler node

D. KSampler Node
The KSampler is the core of the entire workflow, where the entire noise denoising process
occurs, ultimately outputting a latent space image

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Here's an explanation of the KSampler node parameters:


Parameter Name Description Function
model Diffusion model Determines the style and quality of
used for denoising generated images
positive Positive prompt Guides generation to include specified
condition encoding elements
negative Negative prompt Suppresses unwanted content
condition encoding
latent_image Latent space Serves as the input carrier for noise
image to be initialization
denoised
seed Random seed for Controls generation randomness
noise generation
control_after_generate Seed control mode Determines seed variation pattern in batch
after generation generation
steps Number of More steps mean finer details but longer
denoising iterations processing time
cfg Classifier-free Controls prompt constraint strength (too
guidance scale high leads to overfitting)
sampler_name Sampling algorithm Determines the mathematical method for
name denoising path
scheduler Scheduler type Controls noise decay rate and step size
allocation
denoise Denoising strength Controls noise strength added to latent
coefficient space, 0.0 preserves original input
features, 1.0 is complete noise

In the KSampler node, the latent space uses seed as an initialization parameter to construct
random noise, and semantic vectors Positive and Negative are input as conditions to the
diffusion model

Then, based on the number of denoising steps specified by the steps parameter, denoising is
performed. Each denoising step uses the denoising strength coefficient specified by the
denoise parameter to denoise the latent space and generate a new latent space image

E. VAE Decode Node


Converts the latent space image output from the KSampler into a pixel space image

F. Save Image Node

Previews and saves the decoded image from latent space to the local ComfyUI/output folder

Introduction to SD1.5 Model


SD1.5 (Stable Diffusion 1.5) is an AI image generation model developed by Stability AI. It's the
foundational version of the Stable Diffusion series, trained on 512×512 resolution images,
making it particularly good at generating images at this resolution. With a size of about 4GB, it
runs smoothly on consumer-grade GPUs (e.g., 6GB VRAM). Currently, SD1.5 has a rich
ecosystem, supporting various plugins (like ControlNet, LoRA) and optimization tools.
As a milestone model in AI art generation, SD1.5 remains the best entry-level choice thanks to
its open-source nature, lightweight architecture, and rich ecosystem. Although newer versions
like SDXL/SD3 have been released, its value for consumer-grade hardware remains
unmatched.

Basic Information
Release Date: October 2022
Core Architecture: Based on Latent Diffusion Model (LDM)
Training Data: LAION-Aesthetics v2.5 dataset (approximately 590M training steps)
Open Source Features: Fully open-source model/code/training data

Advantages and Limitations


Model Advantages:

Lightweight: Small size, only about 4GB, runs smoothly on consumer GPUs
Low Entry Barrier: Supports a wide range of plugins and optimization tools
Mature Ecosystem: Extensive plugin and tool support
Fast Generation: Smooth operation on consumer GPUs

Model Limitations:

Detail Handling: Hands/complex lighting prone to distortion


Resolution Limits: Quality degrades for direct 1024x1024 generation
Prompt Dependency: Requires precise English descriptions for control

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